Clarification (sorry the question was not specific): They both try to convert the item on the stack to a lua_Number. lua_tonumber will also convert a string that represents a number. How does luaL_checknumber deal with something that's not a number?
There's also luaL_checklong and luaL_checkinteger. Are they the same as (int)luaL_checknumber and (long)luaL_checknumber respectively?
The reference manual does answer this question. I'm citing the Lua 5.2 Reference Manual, but similar text is found in the 5.1 manual as well. The manual is, however, quite terse. It is rare for any single fact to be restated in more than one sentence. Furthermore, you often need to correlate facts stated in widely separated sections to understand the deeper implications of an API function.
This is not a defect, it is by design. This is the reference manual to the language, and as such its primary goal is to completely (and correctly) describe the language.
For more information about "how" and "why" the general advice is to also read Programming in Lua. The online copy is getting rather long in the tooth as it describes Lua 5.0. The current paper edition describes Lua 5.1, and a new edition describing Lua 5.2 is in process. That said, even the first edition is a good resource, as long as you also pay attention to what has changed in the language since version 5.0.
The reference manual has a fair amount to say about the luaL_check* family of functions.
Each API entry's documentation block is accompanied by a token that describes its use of the stack, and under what conditions (if any) it will throw an error. Those tokens are described at section 4.8:
Each function has an indicator like this: [-o, +p, x]
The first field, o, is how many elements the function pops from the
stack. The second field, p, is how many elements the function pushes
onto the stack. (Any function always pushes its results after popping
its arguments.) A field in the form x|y means the function can push
(or pop) x or y elements, depending on the situation; an interrogation
mark '?' means that we cannot know how many elements the function
pops/pushes by looking only at its arguments (e.g., they may depend on
what is on the stack). The third field, x, tells whether the function
may throw errors: '-' means the function never throws any error; 'e'
means the function may throw errors; 'v' means the function may throw
an error on purpose.
At the head of Chapter 5 which documents the auxiliary library as a whole (all functions in the official API whose names begin with luaL_ rather than just lua_) we find this:
Several functions in the auxiliary library are used to check C
function arguments. Because the error message is formatted for
arguments (e.g., "bad argument #1"), you should not use these
functions for other stack values.
Functions called luaL_check* always throw an error if the check is not
satisfied.
The function luaL_checknumber is documented with the token [-0,+0,v] which means that it does not disturb the stack (it pops nothing and pushes nothing) and that it might deliberately throw an error.
The other functions that have more specific numeric types differ primarily in function signature. All are described similarly to luaL_checkint() "Checks whether the function argument arg is a number and returns this number cast to an int", varying the type named in the cast as appropriate.
The function lua_tonumber() is described with the token [-0,+0,-] meaning it has no effect on the stack and does not throw any errors. It is documented to return the numeric value from the specified stack index, or 0 if the stack index does not contain something sufficiently numeric. It is documented to use the more general function lua_tonumberx() which also provides a flag indicating whether it successfully converted a number or not.
It too has siblings named with more specific numeric types that do all the same conversions but cast their results.
Finally, one can also refer to the source code, with the understanding that the manual is describing the language as it is intended to be, while the source is a particular implementation of that language and might have bugs, or might reveal implementation details that are subject to change in future versions.
The source to luaL_checknumber() is in lauxlib.c. It can be seen to be implemented in terms of lua_tonumberx() and the internal function tagerror() which calls typerror() which is implemented with luaL_argerror() to actually throw the formatted error message.
They both try to convert the item on the stack to a lua_Number. lua_tonumber will also convert a string that represents a number. luaL_checknumber throws a (Lua) error when it fails a conversion - it long jumps and never returns from the POV of the C function. lua_tonumber merely returns 0 (which can be a valid return as well.) So you could write this code which should be faster than checking with lua_isnumber first.
double r = lua_tonumber(_L, idx);
if (r == 0 && !lua_isnumber(_L, idx))
{
// Error handling code
}
return r;
Related
I was studying one llvm code, where i have found a line,
const ExtractValueInst *EI = cast<ExtractValueInst>(I);
st.setValue(I, st.getValue(EI->getAggregateOperand()));
Now, I understand why cast<...> is being used, but I can't relate with the ExtarctValueInst, Can you give me one example what is this instruction in IR and what is the equivalence C code? and also I want to know about getAggregateOperand() function also. Thank you in advance.
Suppose you have a function that returns a 32-bit integer that isn't really one integer, but rather a set of smaller bitfields and/or bools. Perhaps the bottom six bits are an integer in the range 0-63, the seventh bit is a boolean, etc.
Somewhere you have a call to that function, and a little below the call, you have code that uses the i6 that's a part of the return value. So you create an extractvalue to extract a value from the composite return value. (If the composite were in main memory you'd probably create a getelementptr/load pair, but since it's most likely in a CPU register you create an extractvalue.)
This is quite often used e.g. in exception handling; a catch clause has a single parameter which is a composite of two things, and the catch clause tests one of the two components to determine whether to catch the exception or pass it on.
increment([]) -> [];
increment([H|T]) -> [H+1|increment(T)].
decrement([]) -> [];
decrement([H|T]) -> [H-1|decrement(T)].
So I have this code but I don't know how they properly work like in java.
Java and Erlang are different beasts. I don't recommend trying to make comparisons to Java when learning Erlang, especially if Java is the only language you know so far. The code you've posted is a good example of the paradigm known as "functional programming". I'd suggest doing some reading on that subject to help you understand what's going on. To try to break this down as far as Erlang goes, you need to understand that an Erlang function is completely different from a Java method.
In Java, your method signature is composed of the method name and the types of its arguments. The return type can also be significant. A Java increment method like the function you wrote might be written like List<Integer> increment(List<Integer> input). The body of the Java method would probably iterate through the list an element at a time and set each element to itself plus one:
List<Integer> increment(List<Integer> input) {
for (int i = 0; i < input.size; i++) {
input.set(i, input.get(i) + 1);
}
}
Erlang has almost nothing in common with this. To begin with, an erlang function's "signature" is the name and arity of the function. Arity means how many arguments the function accepts. So your increment function is known as increment/1, and that's its unique signature. The way you write the argument list inside the parentheses after the function name has less to do with argument types than with the pattern of the data passed to it. A function like increment([]) -> ... can only successfully be called by passing it [], the empty list. Likewise, the function increment([Item]) -> ... can only be successfully called by passing it a list with one item in it, and increment([Item1, Item2]) -> ... must be passed a list with two items in it. This concept of matching data to patterns is quite aptly known as "pattern matching", and you'll find it in many functional languages. In Erlang functions, it's used to select which head of the function to execute. This bears a rough similarity to Java's method overloading, where you can have many methods with the same name but different argument types; however a pattern in an Erlang function head can bind variables to different pieces of the arguments that match the pattern.
In your code example, the function increment/1 has two heads. The first head is executed only if you pass an empty list to the function. The second head is executed only if you pass a non-empty list to the function. When that happens, two variables, H and T, are bound. H is bound to the first item of the list, and T is bound to the rest of the list, meaning all but the first item. That's because the pattern [H|T] matches a non-empty list, including a list with one element, in which case T would be bound to the empty list. The variables thus bound can be used in the body of the function.
The bodies of your functions are a very typical form of iterating a list in Erlang to produce a new list. It's typical because of another important difference from Java, which is that Erlang data is immutable. That means there's no such concept as "setting an element of a list" like I did in the Java code above. If you want to change a list, you have to build a new one, which is what your code does. It effectively says:
The result of incrementing the empty list is the empty list.
The result of incrementing a non-empty list is:
Take the first element of the list: H.
Increment the rest of the list: increment(T).
Prepend H+1 to the result of incrementing the rest of the list.
Note that you want to be careful about how you build lists in Erlang, or you can end up wasting a lot of resources. The List Handling User's Guide is a good place to learn about that. Also note that this code uses a concept known as "recursion", meaning that the function calls itself. In many popular languages, including Java, recursion is of limited usefulness because each new function call adds a stack frame, and your available memory space for stack frames is relatively limited. Erlang and many functional languages support a thing known as "tail call elimination", which is a feature that allows properly written code to recurse indefinitely without exhausting any resources.
Hopefully this helps explain things. If you can ask a more specific question, you might get a better answer.
I am starting to work with DLV (Disjunctive Datalog) and I have a rule that is reporting a "Rule is not safe" error, when running the code. The rule is the following:
foo(R, 1) :- not foo(R, _)
I have read the manual and seen that "cyclic dependencies are disallowed". I guess this is why I am being reported the error, but I am not sure how this statement is so problematic to DLV. The final goal is to have some kind of initialization in case that the predicate has not been defined.
More precisely, if there is no occurrence of 'foo' with the parameter R (and anything else), then define it with parameters R and 1. Once it is defined, the rule shouldn't be triggered again. So, this is not a real recursion in my opinion.
Any comments on how to solve this issue are welcomed!
I have realised that I probably need another predicate to match the parameter R in the body of the rule. Something like this:
foo(R, 1) :- not foo(R, _), bar(R)
Since, otherwise there would be no way to know whether there are no occurrences of foo(R, _). I don't know whether I made myself clear.
Anyway, this doesn't work either :(
To the particular "Rule is not safe" error: First of all this has nothing to do with cyclic or acyclic dependencies. The same error message shows up for the non-cyclic program:
foo2(R, 1) :- not foo(R,_), bar(R).
The problem is that the program is actually not safe (http://www.dlvsystem.com/html/DLV_User_Manual.html#SAFETY). As mentioned in the section on negative rules (anchor #AEN375, I am only allowed to use 2 links in my answer):
Variables, which occur in a negated literal, must also occur in a
positive literal in the body.
Observe that the _ is an anonymous variable. I.e., the program
foo(R,1) :- not foo(R,_), bar(R).
can be equivalently written as (and is equivalent to)
foo(R,1) :- not foo(R,X), bar(R).
Anonymous variables (DLV manual, anchor #AEN264 - at the end of the section) just allow us to avoid inventing names for variables that will only occur once within the rule (i.e. for variables that only express "there is some value, I absolutely do not care about it), but they are variables nevertheless. And since negation with not is "negation" and not "true negation" (or "strong negation" as it is also often called), none of the three safety conditions is satisfied by the rule.
A very rough and high-level intuition for safety is that it guarantees that every variable in the program can be assigned to some finite domain - as it is now the case with R by adding bar(R). However, the same also must be the case for the anonymous variable _ .
To the actual problem of defining default values:
As pointed out by lambda.xy.x, the problem here is the Answer Set (or stable model) semantics of DLV: Trying to do it in one rule does not give any solution:
In order to get a safe program, we could replace the above problems e.g. by
foo(1,2). bar(1). bar(2).
tmp(R) :- foo(R,_).
foo(R,1) :- not tmp(R), bar(R).
This has no stable model:
Assume the answer is, as intended,
{foo(1,2), bar(1), bar(2), foo(2,1)}
However, this is not a valid model, since tmp(R) :- foo(R,_) would require it to contain tmp(2). But then, "not tmp(2)" is no longer true, and therefore having foo(2,1) in the model violates the required minimality of the model. (This is not exactly what is going on, more a rough intuition. More technical details could be found in any article on answer set programming, a quick Google search gave me this paper as one of the first results: http://www.kr.tuwien.ac.at/staff/tkren/pub/2009/rw2009-asp.pdf)
In order to solve the problem, it is therefore somehow necessary to "break the cycle". One possibility would be:
foo(1,2). bar(1). bar(2). bar(3).
tmp(R) :- foo(R,X), X!=1.
foo(R,1) :- bar(R), not tmp(R).
I.e., by explicitly stating that we want to add R into the intermediate atom only if the value is different from 1, having foo(2,1) in the model does not contradict tmp(2) not being part of the model as well. Of course, this no longer allows to distinguish whether foo(R,1) is there as default value or by input, but if this is not required ...
Another possibility would be to not use foo for the computation, but some foo1 instead. I.e. having
foo1(R,X) :- foo(R,X).
tmp(R) :- foo(R,_).
foo1(R,1) :- bar(R), not tmp(R).
and then just use foo1 instead of foo.
I'm reading through two articles right now and am a little confused.
This article - http://blog.golang.org/laws-of-reflection says
> var r io.Reader
tty, err := os.OpenFile("/dev/tty", os.O_RDWR, 0)
if err != nil {
return nil, err
}
r = tty
r contains, schematically, the (value, type) pair, (tty, *os.File).
Notice that the type *os.File implements methods other than Read; even
though the interface value provides access only to the Read method,
the value inside carries all the type information about that value.
This other article, says
In terms of our example, the itable for Stringer holding type Binary
lists the methods used to satisfy Stringer, which is just String:
Binary's other methods (Get) make no appearance in the itable.
It feels like these two are in opposition. According to the second article, the variable r in the first extract should be (tty, io.Reader), as that is the static type of r. Instead, the article says that *os.File is the type of tty. If the second example were right, then the diagram in the first example should have all of the methods implemented by the Binary type.
Where am I going wrong?
The two articles are explaining a similar concept at two very different levels of granularity. As Volker said, "Laws of Reflection" is a basic overview of what is actually happening when you examine objects via reflection. Your second article is examining the dynamic dispatch properties of an interface (which, can be resolved via reflection as well) and how the runtime resolves them at runtime.
According to the second article, the variable r in the first extract should be (tty, io.Reader)
Given that understanding, think of an interface at runtime as a "wrapper object". It exists to provide information about another object (the itable from your second article) to know where to jump to in the wrapped objects layout (implementation may differ between versions .. but the principle is basically the same for most languages).
This is why calling Read on r works .. first it will check the itable and jump to the function that is laid out for the os.File type. If that was an interface .. you would be looking at another dereference and dispatch (which IIRC isn't applicable at all in Go).
RE: Reflection - you're getting an easily digestible representation of this, in the form of a (value, type) pair (via the reflect.ValueOf and reflect.TypeOf methods).
Both are correct, r "holds" (tty, *os.File) and that is what the second article says. Note that Laws of Reflection is a bit more high-level and does not mention implementation details (which might change in each release) like discussed in the second article. The diagram of the second article read as follows: "s contains schematically (b, *Binary). s is of type Stringer, its data is a Binary with value 200 and s's itable contains one method String and the other methods of Binary (or *Binary) are not represented in the itable and thus not accessible for s.
Note that I think the actual implementation of interfaces in Go 1.4 is different from what the second article (Russ's?) states. "
if I call a function foo(t[1]) through the C API, can I in any way see what table and what index is as argument, in this case t and 1?
The problem at hand is a function move_card(card, table_slots[0]) where I move a card from one slot on a game area to another one. table_slots can also be hand or player_slots. This can be solved using a metatable, stating the name of the table being accessed. But the index is impossible to solve, or is it? table_slots[0] can be a card table, or an overlay or stack (arrays of cards), or nil if it's empty, as could player_slots. But what I need to know is in fact if it's from a table_slots or player_slots.
Could I hack some code analysis? Like, get the line where the function call is made, and then grep the index through a regexp? I could send table_slots[0] as a string, also, like move_card(card, "table_slot[0]"). Not as elegant, but still working.
No, once a value reaches a function, be it C or Lua, its origins are lost.
On the other hand, when an error occurs in Lua, the runtime system tries quite a bit to reconstruct the origin of the relevant values, but it does not always succeeds.